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Diffusers

diffusers

Solving everything with diffusion models!

Models in Library1
Datasets in Library0
Models on Hugging Face81
Followers722

Models

basemodel: stabilityai/stable-diffusion-xl-base-1.0 - stable-diffusion-xl - stable-diffusion-xl-diffusers - text-to-image - diffusers - inpainting SD-XL Inpainting 0.1 is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask. The SD-XL Inpainting 0.1 was initialized with the stable-diffusion-xl-base-1.0 weights. The model is trained for 40k steps at resolution 1024x1024 and 5% dropping of the text-conditioning to improve classifier-free classifier-free guidance sampling. For inpainting, the UNet has 5 additional input channels (4 for the encoded masked-image and 1 for the…

Open weights openrail++ 2.6B parameters diffusers